AI projects from Hugging Face

OpenDDE

aurekaresearch · Kit

Predicts the 3D structure of proteins and other biological molecules down to each atom, for early-stage drug research. Runs from a command line on a GPU, a CPU…

Tessera-small

thejagstudio · Kit

Generates 16×16 pixel art sprites from a text prompt with a small diffusion model that fills in all 256 pixels in about 16 steps. For non-commercial game and a…

MOSS-SoundEffect-v2.0

OpenMOSS-Team · Kit

Generates sound effects from a text prompt, such as nature, city scenes, animals and human actions, up to 30 seconds at 48 kHz. For creators and developers who…

K2-Horizon-7B-Uno

IFM · Kit

A LoRA adapter that adds a diffusion pathway to IFM's K2-Horizon-7B language model, for research into faster text generation. It holds only the adapter, not th…

ACE-Step 1.5 XL Turbo

ACE-Step · Kit

Generates music and other audio in just 8 steps with a 4B-parameter diffusion model, run locally through a Gradio interface. For creators and developers with a…

sd-vae-ft-mse-original

stabilityai · Kit

Replaces the autoencoder in Stable Diffusion with a fine-tuned decoder that reconstructs images, especially faces, more smoothly. Suits people running the orig…

FlashVSR+

ruchangyao · Kit

A one-step diffusion model that sharpens low-quality video into higher-resolution video, with a fast Turbo variant and a higher-quality Base variant. For video…

FLUX.2 Klein 4B for ComfyUI

Comfy-Org · Kit

Repackaged FLUX.2 Klein 4B model files for ComfyUI that turn a text prompt into an image. Includes diffusion models, text encoders and a VAE, with a ready text…

ACE-Step 1.5 by Comfy-Org

Comfy-Org · Kit

Model files for generating music from text in ComfyUI, repackaged from the original ACE-Step 1.5 release. Includes checkpoint, diffusion, text encoder and VAE…

MOSS-Image-SFT-GGUF

unsloth · Kit

Runs the MOSS-Image text to image model from compressed GGUF files on CPU, Mac or GPU, through Unsloth Studio or stable-diffusion.cpp. For people who make imag…

Iris-3B by speridlabs

speridlabs · Kit

Generates 1024-pixel images from text prompts by painting every pixel directly, with no VAE or latent space. Fine-tuned versions estimate depth and restore or…

Animagine XL 4.0

cagliostrolab · Kit

Generates anime-style images from tag-based text prompts, a Stable Diffusion XL fine-tune retrained on a large anime dataset. For artists and developers using…

Juggernaut-XL-v9

RunDiffusion · Kit

Makes photorealistic images from text prompts, as a Stable Diffusion XL fine-tune that the card says runs on 8 GB of video memory. For creators who want photo-…

ACE-Step 1.5

ACE-Step · Kit

Generates full songs from a text description, with a language model that plans lyrics and structure and a diffusion transformer that renders the audio. Built t…

stable-diffusion-xl-base-1.0

stabilityai · Kit

Stability AI's SDXL base model, which makes and edits images from text prompts through the diffusers library, alone or paired with a refiner.

HairDeck

hyz317 · Kit

Turns one character image into stylized 3D hair cards with colors and spring-bone chains for animation, exported as GLB meshes or a VRM avatar. For avatar make…

VoxCPM2

openbmb · Kit

Speaks text in 30 languages and nine Chinese dialects at 48 kHz, designs a new voice from a written description and clones a voice from a short clip. An open 2…

Diffusion-Qwen3.8-27B

summerAIO · Kit

Generates text by parallel unmasking from Qwen3.8-27B converted into a masked diffusion language model. An experimental research checkpoint for researchers and…

pjit-21k

Linum-AI · Kit

Linum's experimental 2B text to image model, trained for about $21K of compute, that works best with dense 50 to 80 word prompts and cannot write text in image…

MuseTalkDemo

TMElyralab · Kit

A Docker-packaged demo app from TMElyralab for its MuseTalk model, bundled with DWPose, a Stable Diffusion autoencoder and ffmpeg tooling. It suits developers…