Mindket

Kits of AI agent skills, rules and prompts

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Nanosaur2-670M

well9472 · Kit

Makes illustrations from text prompts with a small 670M-parameter diffusion model, run through ComfyUI. Built as a low-budget research model for people who wan…

New

Stable Diffusion v1-5

stable-diffusion-v1-5 · Kit

Generates images from text prompts with a latent diffusion model of 860M parameters, usable in diffusers, ComfyUI and other tools. Suits researchers and develo…

New

stable-diffusion-inpainting

stable-diffusion-v1-5 · Kit

Generates photo-realistic images from text and fills masked areas of a picture, using Stable Diffusion with an extra mask input, in diffusers or AUTOMATIC1111.…

New

Juggernaut XI v11

RunDiffusion · Kit

A text-to-image model from RunDiffusion, part of the Juggernaut SDXL line, tuned for photo-style pictures with better hands, faces and prompt following. For im…

New

Primal VAE

ostris · Kit

Provides small image autoencoders that compress pictures to one or two latent channels and need no scaling factor, then patch up to match the Flux or Stable Di…

New

RealVisXL V5.0 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of RealVisXL V5.0 with stable-diffusion.cpp. For people who want local image generation wit…

New

Juggernaut XL v9 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of Juggernaut XL v9 with stable-diffusion.cpp. For people who want local image generation w…

New

LongCat-AudioDiT-3.5B

meituan-longcat · Kit

Turns Chinese or English text into speech with a diffusion model that works directly on the waveform latent space, and can clone a voice from a reference sampl…

New

Steers Stable Diffusion 1.5 with a depth map so generated images keep the layout of a source picture. ControlNet v1.1 in diffusers format, for artists and deve…

New

Generates images from text on modest GPUs, using quantized Z-Image-Turbo files made for stable-diffusion.cpp. For people who run image models on their own comp…

New

DiT-Wikiart-Small

kaupane · Kit

Predicts noise for a small diffusion transformer trained from scratch on Wikiart paintings, built to make art images from a genre and a style. For developers a…

New

FlashVSR

JunhaoZhuang · Kit

A one-step diffusion model that upscales low-resolution video as a stream, at about 17 frames per second for 768 x 1408 video on one A100 GPU. This is the firs…

New

InstantID by InstantX

InstantX · Kit

Generates new images of a person from one face photo and a text prompt, keeping the face recognizable. Made for developers and artists who run Stable Diffusion…

New

Archived Stable Diffusion v1.5 checkpoints for ComfyUI that turn a text prompt into an image, kept for legacy testing. Includes the original file and an FP16 c…

New

Cola DLM

ByteDance-Seed · Kit

A research text-generation model that writes through a continuous latent space, pairing a text autoencoder with a diffusion transformer. For researchers and de…

New

MOSS-Image-SFT GGUF

axiomofmind · Kit

Unofficial GGUF files of the MOSS-Image-SFT diffusion transformer in Q8_0 and Q4_K_S, for running text-to-image locally with stable-diffusion.cpp. Needs a sepa…

New

Tessera-small

thejagstudio · Kit

Generates 16×16 pixel art sprites from a text prompt with a small diffusion model that fills in all 256 pixels in about 16 steps. For non-commercial game and a…

New

A LoRA adapter that adds a diffusion pathway to IFM's K2-Horizon-7B language model, for research into faster text generation. It holds only the adapter, not th…

New

ACE-Step 1.5 XL Turbo

ACE-Step · Kit

Generates music and other audio in just 8 steps with a 4B-parameter diffusion model, run locally through a Gradio interface. For creators and developers with a…

New

sd-vae-ft-mse-original

stabilityai · Kit

Replaces the autoencoder in Stable Diffusion with a fine-tuned decoder that reconstructs images, especially faces, more smoothly. Suits people running the orig…

New