AI skills and agents for developers

Stable Diffusion v1-5

stable-diffusion-v1-5 · Kit

Generates images from text prompts with a latent diffusion model of 860M parameters, usable in diffusers, ComfyUI and other tools. Suits researchers and develo…

stable-diffusion-inpainting

stable-diffusion-v1-5 · Kit

Generates photo-realistic images from text and fills masked areas of a picture, using Stable Diffusion with an extra mask input, in diffusers or AUTOMATIC1111.…

Juggernaut XI v11

RunDiffusion · Kit

A text-to-image model from RunDiffusion, part of the Juggernaut SDXL line, tuned for photo-style pictures with better hands, faces and prompt following. For im…

Primal VAE

ostris · Kit

Provides small image autoencoders that compress pictures to one or two latent channels and need no scaling factor, then patch up to match the Flux or Stable Di…

RealVisXL V5.0 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of RealVisXL V5.0 with stable-diffusion.cpp. For people who want local image generation wit…

Juggernaut XL v9 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of Juggernaut XL v9 with stable-diffusion.cpp. For people who want local image generation w…

LongCat-AudioDiT-3.5B

meituan-longcat · Kit

Turns Chinese or English text into speech with a diffusion model that works directly on the waveform latent space, and can clone a voice from a reference sampl…

ControlNet v1.1 depth (SD 1.5)

lllyasviel · Kit

Steers Stable Diffusion 1.5 with a depth map so generated images keep the layout of a source picture. ControlNet v1.1 in diffusers format, for artists and deve…

DiT-Wikiart-Small

kaupane · Kit

Predicts noise for a small diffusion transformer trained from scratch on Wikiart paintings, built to make art images from a genre and a style. For developers a…

FlashVSR

JunhaoZhuang · Kit

A one-step diffusion model that upscales low-resolution video as a stream, at about 17 frames per second for 768 x 1408 video on one A100 GPU. This is the firs…

InstantID by InstantX

InstantX · Kit

Generates new images of a person from one face photo and a text prompt, keeping the face recognizable. Made for developers and artists who run Stable Diffusion…

Stable Diffusion v1.5 Archive

Comfy-Org · Kit

Archived Stable Diffusion v1.5 checkpoints for ComfyUI that turn a text prompt into an image, kept for legacy testing. Includes the original file and an FP16 c…

MOSS-Image-SFT GGUF

axiomofmind · Kit

Unofficial GGUF files of the MOSS-Image-SFT diffusion transformer in Q8_0 and Q4_K_S, for running text-to-image locally with stable-diffusion.cpp. Needs a sepa…

Tessera-small

thejagstudio · Kit

Generates 16×16 pixel art sprites from a text prompt with a small diffusion model that fills in all 256 pixels in about 16 steps. For non-commercial game and a…

K2-Horizon-7B-Uno

IFM · Kit

A LoRA adapter that adds a diffusion pathway to IFM's K2-Horizon-7B language model, for research into faster text generation. It holds only the adapter, not th…

sd-vae-ft-mse-original

stabilityai · Kit

Replaces the autoencoder in Stable Diffusion with a fine-tuned decoder that reconstructs images, especially faces, more smoothly. Suits people running the orig…

FlashVSR+

ruchangyao · Kit

A one-step diffusion model that sharpens low-quality video into higher-resolution video, with a fast Turbo variant and a higher-quality Base variant. For video…

FLUX.2 Klein 4B for ComfyUI

Comfy-Org · Kit

Repackaged FLUX.2 Klein 4B model files for ComfyUI that turn a text prompt into an image. Includes diffusion models, text encoders and a VAE, with a ready text…

ACE-Step 1.5 by Comfy-Org

Comfy-Org · Kit

Model files for generating music from text in ComfyUI, repackaged from the original ACE-Step 1.5 release. Includes checkpoint, diffusion, text encoder and VAE…

Animagine XL 4.0

cagliostrolab · Kit

Generates anime-style images from tag-based text prompts, a Stable Diffusion XL fine-tune retrained on a large anime dataset. For artists and developers using…