Mindket

AI projects from Hugging Face

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Stable Diffusion v1-5

stable-diffusion-v1-5 · Kit

Generates images from text prompts with a latent diffusion model of 860M parameters, usable in diffusers, ComfyUI and other tools. Suits researchers and develo…

New

stable-diffusion-inpainting

stable-diffusion-v1-5 · Kit

Generates photo-realistic images from text and fills masked areas of a picture, using Stable Diffusion with an extra mask input, in diffusers or AUTOMATIC1111.…

New

Juggernaut XI v11

RunDiffusion · Kit

A text-to-image model from RunDiffusion, part of the Juggernaut SDXL line, tuned for photo-style pictures with better hands, faces and prompt following. For im…

New

RealVisXL V5.0 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of RealVisXL V5.0 with stable-diffusion.cpp. For people who want local image generation wit…

New

Juggernaut XL v9 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of Juggernaut XL v9 with stable-diffusion.cpp. For people who want local image generation w…

New

Steers Stable Diffusion 1.5 with a depth map so generated images keep the layout of a source picture. ControlNet v1.1 in diffusers format, for artists and deve…

New

Generates images from text on modest GPUs, using quantized Z-Image-Turbo files made for stable-diffusion.cpp. For people who run image models on their own comp…

New

MOSS-Image-SFT GGUF

axiomofmind · Kit

Unofficial GGUF files of the MOSS-Image-SFT diffusion transformer in Q8_0 and Q4_K_S, for running text-to-image locally with stable-diffusion.cpp. Needs a sepa…

New

sd-vae-ft-mse-original

stabilityai · Kit

Replaces the autoencoder in Stable Diffusion with a fine-tuned decoder that reconstructs images, especially faces, more smoothly. Suits people running the orig…

New

MOSS-Image-SFT-GGUF

unsloth · Kit

Runs the MOSS-Image text to image model from compressed GGUF files on CPU, Mac or GPU, through Unsloth Studio or stable-diffusion.cpp. For people who make imag…

New

Animagine XL 4.0

cagliostrolab · Kit

Generates anime-style images from tag-based text prompts, a Stable Diffusion XL fine-tune retrained on a large anime dataset. For artists and developers using…

New

Juggernaut-XL-v9

RunDiffusion · Kit

Makes photorealistic images from text prompts, as a Stable Diffusion XL fine-tune that the card says runs on 8 GB of video memory. For creators who want photo-…

New

Stability AI's SDXL base model, which makes and edits images from text prompts through the diffusers library, alone or paired with a refiner.

New