Stable Diffusion v1-5
stable-diffusion-v1-5 · Kit
Generates images from text prompts with a latent diffusion model of 860M parameters, usable in diffusers, ComfyUI and other tools. Suits researchers and develo…
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AI skills and agents for developers
Prompts, coding agent briefs and kits for writing and reviewing code with fewer tokens.
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stable-diffusion-inpainting
stable-diffusion-v1-5 · Kit
Generates photo-realistic images from text and fills masked areas of a picture, using Stable Diffusion with an extra mask input, in diffusers or AUTOMATIC1111.…
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Juggernaut XI v11
RunDiffusion · Kit
A text-to-image model from RunDiffusion, part of the Juggernaut SDXL line, tuned for photo-style pictures with better hands, faces and prompt following. For im…
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Primal VAE
ostris · Kit
Provides small image autoencoders that compress pictures to one or two latent channels and need no scaling factor, then patch up to match the Flux or Stable Di…
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RealVisXL V5.0 GGUF
off-grid-ai · Kit
Generates pictures from text on your own computer using GGUF copies of RealVisXL V5.0 with stable-diffusion.cpp. For people who want local image generation wit…
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Juggernaut XL v9 GGUF
off-grid-ai · Kit
Generates pictures from text on your own computer using GGUF copies of Juggernaut XL v9 with stable-diffusion.cpp. For people who want local image generation w…
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LongCat-AudioDiT-3.5B
meituan-longcat · Kit
Turns Chinese or English text into speech with a diffusion model that works directly on the waveform latent space, and can clone a voice from a reference sampl…
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ControlNet v1.1 depth (SD 1.5)
lllyasviel · Kit
Steers Stable Diffusion 1.5 with a depth map so generated images keep the layout of a source picture. ControlNet v1.1 in diffusers format, for artists and deve…
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DiT-Wikiart-Small
kaupane · Kit
Predicts noise for a small diffusion transformer trained from scratch on Wikiart paintings, built to make art images from a genre and a style. For developers a…
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FlashVSR
JunhaoZhuang · Kit
A one-step diffusion model that upscales low-resolution video as a stream, at about 17 frames per second for 768 x 1408 video on one A100 GPU. This is the firs…
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InstantID by InstantX
InstantX · Kit
Generates new images of a person from one face photo and a text prompt, keeping the face recognizable. Made for developers and artists who run Stable Diffusion…
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Stable Diffusion v1.5 Archive
Comfy-Org · Kit
Archived Stable Diffusion v1.5 checkpoints for ComfyUI that turn a text prompt into an image, kept for legacy testing. Includes the original file and an FP16 c…
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MOSS-Image-SFT GGUF
axiomofmind · Kit
Unofficial GGUF files of the MOSS-Image-SFT diffusion transformer in Q8_0 and Q4_K_S, for running text-to-image locally with stable-diffusion.cpp. Needs a sepa…
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Tessera-small
thejagstudio · Kit
Generates 16×16 pixel art sprites from a text prompt with a small diffusion model that fills in all 256 pixels in about 16 steps. For non-commercial game and a…
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K2-Horizon-7B-Uno
IFM · Kit
A LoRA adapter that adds a diffusion pathway to IFM's K2-Horizon-7B language model, for research into faster text generation. It holds only the adapter, not th…
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sd-vae-ft-mse-original
stabilityai · Kit
Replaces the autoencoder in Stable Diffusion with a fine-tuned decoder that reconstructs images, especially faces, more smoothly. Suits people running the orig…
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FlashVSR+
ruchangyao · Kit
A one-step diffusion model that sharpens low-quality video into higher-resolution video, with a fast Turbo variant and a higher-quality Base variant. For video…
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FLUX.2 Klein 4B for ComfyUI
Comfy-Org · Kit
Repackaged FLUX.2 Klein 4B model files for ComfyUI that turn a text prompt into an image. Includes diffusion models, text encoders and a VAE, with a ready text…
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ACE-Step 1.5 by Comfy-Org
Comfy-Org · Kit
Model files for generating music from text in ComfyUI, repackaged from the original ACE-Step 1.5 release. Includes checkpoint, diffusion, text encoder and VAE…
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Animagine XL 4.0
cagliostrolab · Kit
Generates anime-style images from tag-based text prompts, a Stable Diffusion XL fine-tune retrained on a large anime dataset. For artists and developers using…
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ACE-Step 1.5
ACE-Step · Kit
Generates full songs from a text description, with a language model that plans lyrics and structure and a diffusion transformer that renders the audio. Built t…
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stable-diffusion-xl-base-1.0
stabilityai · Kit
Stability AI's SDXL base model, which makes and edits images from text prompts through the diffusers library, alone or paired with a refiner.
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MuseTalkDemo
TMElyralab · Kit
A Docker-packaged demo app from TMElyralab for its MuseTalk model, bundled with DWPose, a Stable Diffusion autoencoder and ffmpeg tooling. It suits developers…
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LGM-Mini
dylanebert · Kit
Runs a small Gradio demo built on the LGM model and several multi-view diffusion models, on a Hugging Face GPU. It suits people interested in 3D generation res…
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Texture-Craft
ArshPathan · Kit
Makes transparent 16x16 item textures from text prompts with a diffusion model, through a Gradio Space, a local REST server or an MCP tool. For game makers and…
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InstantID
InstantX · Kit
A Gradio demo built around the InstantX/InstantID model, with a Stable Diffusion XL pipeline, ControlNet depth and canny models and style templates, online or…
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P-JiT 21K
hugging-apps · Kit
Generates images from text prompts with P-JiT 21K, a diffusion transformer that works directly on pixels. Includes a prompt expander for short prompts. For ima…
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