AI projects from Hugging Face

Nanosaur2-670M

well9472 · Kit

Makes illustrations from text prompts with a small 670M-parameter diffusion model, run through ComfyUI. Built as a low-budget research model for people who wan…

AuK

tencent · Kit

A 1.5B speech model from Tencent that generates speech from text, edits recordings, and also cleans up audio and separates sources, all through plain-language…

Stable Diffusion v1-5

stable-diffusion-v1-5 · Kit

Generates images from text prompts with a latent diffusion model of 860M parameters, usable in diffusers, ComfyUI and other tools. Suits researchers and develo…

stable-diffusion-inpainting

stable-diffusion-v1-5 · Kit

Generates photo-realistic images from text and fills masked areas of a picture, using Stable Diffusion with an extra mask input, in diffusers or AUTOMATIC1111.…

Juggernaut XI v11

RunDiffusion · Kit

A text-to-image model from RunDiffusion, part of the Juggernaut SDXL line, tuned for photo-style pictures with better hands, faces and prompt following. For im…

Primal VAE

ostris · Kit

Provides small image autoencoders that compress pictures to one or two latent channels and need no scaling factor, then patch up to match the Flux or Stable Di…

RealVisXL V5.0 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of RealVisXL V5.0 with stable-diffusion.cpp. For people who want local image generation wit…

Juggernaut XL v9 GGUF

off-grid-ai · Kit

Generates pictures from text on your own computer using GGUF copies of Juggernaut XL v9 with stable-diffusion.cpp. For people who want local image generation w…

LongCat-AudioDiT-3.5B

meituan-longcat · Kit

Turns Chinese or English text into speech with a diffusion model that works directly on the waveform latent space, and can clone a voice from a reference sampl…

ControlNet v1.1 depth (SD 1.5)

lllyasviel · Kit

Steers Stable Diffusion 1.5 with a depth map so generated images keep the layout of a source picture. ControlNet v1.1 in diffusers format, for artists and deve…

Z-Image-Turbo-GGUF by leejet

leejet · Kit

Generates images from text on modest GPUs, using quantized Z-Image-Turbo files made for stable-diffusion.cpp. For people who run image models on their own comp…

DiT-Wikiart-Small

kaupane · Kit

Predicts noise for a small diffusion transformer trained from scratch on Wikiart paintings, built to make art images from a genre and a style. For developers a…

FlashVSR

JunhaoZhuang · Kit

A one-step diffusion model that upscales low-resolution video as a stream, at about 17 frames per second for 768 x 1408 video on one A100 GPU. This is the firs…

InstantID by InstantX

InstantX · Kit

Generates new images of a person from one face photo and a text prompt, keeping the face recognizable. Made for developers and artists who run Stable Diffusion…

Stable Diffusion v1.5 Archive

Comfy-Org · Kit

Archived Stable Diffusion v1.5 checkpoints for ComfyUI that turn a text prompt into an image, kept for legacy testing. Includes the original file and an FP16 c…

GenFirst

ByteDance-Seed · Kit

Research weights from ByteDance-Seed for image generators whose VAE and prior are trained together. They cover class-conditional and text-to-image sampling, fo…

Cola DLM

ByteDance-Seed · Kit

A research text-generation model that writes through a continuous latent space, pairing a text autoencoder with a diffusion transformer. For researchers and de…

MOSS-Image-SFT GGUF

axiomofmind · Kit

Unofficial GGUF files of the MOSS-Image-SFT diffusion transformer in Q8_0 and Q4_K_S, for running text-to-image locally with stable-diffusion.cpp. Needs a sepa…

OpenDDE

aurekaresearch · Kit

Predicts the 3D structure of proteins and other biological molecules down to each atom, for early-stage drug research. Runs from a command line on a GPU, a CPU…

Tessera-small

thejagstudio · Kit

Generates 16×16 pixel art sprites from a text prompt with a small diffusion model that fills in all 256 pixels in about 16 steps. For non-commercial game and a…

MOSS-SoundEffect-v2.0

OpenMOSS-Team · Kit

Generates sound effects from a text prompt, such as nature, city scenes, animals and human actions, up to 30 seconds at 48 kHz. For creators and developers who…

K2-Horizon-7B-Uno

IFM · Kit

A LoRA adapter that adds a diffusion pathway to IFM's K2-Horizon-7B language model, for research into faster text generation. It holds only the adapter, not th…

ACE-Step 1.5 XL Turbo

ACE-Step · Kit

Generates music and other audio in just 8 steps with a 4B-parameter diffusion model, run locally through a Gradio interface. For creators and developers with a…

sd-vae-ft-mse-original

stabilityai · Kit

Replaces the autoencoder in Stable Diffusion with a fine-tuned decoder that reconstructs images, especially faces, more smoothly. Suits people running the orig…

FlashVSR+

ruchangyao · Kit

A one-step diffusion model that sharpens low-quality video into higher-resolution video, with a fast Turbo variant and a higher-quality Base variant. For video…

FLUX.2 Klein 4B for ComfyUI

Comfy-Org · Kit

Repackaged FLUX.2 Klein 4B model files for ComfyUI that turn a text prompt into an image. Includes diffusion models, text encoders and a VAE, with a ready text…

ACE-Step 1.5 by Comfy-Org

Comfy-Org · Kit

Model files for generating music from text in ComfyUI, repackaged from the original ACE-Step 1.5 release. Includes checkpoint, diffusion, text encoder and VAE…

MOSS-Image-SFT-GGUF

unsloth · Kit

Runs the MOSS-Image text to image model from compressed GGUF files on CPU, Mac or GPU, through Unsloth Studio or stable-diffusion.cpp. For people who make imag…

Iris-3B by speridlabs

speridlabs · Kit

Generates 1024-pixel images from text prompts by painting every pixel directly, with no VAE or latent space. Fine-tuned versions estimate depth and restore or…

Animagine XL 4.0

cagliostrolab · Kit

Generates anime-style images from tag-based text prompts, a Stable Diffusion XL fine-tune retrained on a large anime dataset. For artists and developers using…

Juggernaut-XL-v9

RunDiffusion · Kit

Makes photorealistic images from text prompts, as a Stable Diffusion XL fine-tune that the card says runs on 8 GB of video memory. For creators who want photo-…

ACE-Step 1.5

ACE-Step · Kit

Generates full songs from a text description, with a language model that plans lyrics and structure and a diffusion transformer that renders the audio. Built t…

stable-diffusion-xl-base-1.0

stabilityai · Kit

Stability AI's SDXL base model, which makes and edits images from text prompts through the diffusers library, alone or paired with a refiner.

HairDeck

hyz317 · Kit

Turns one character image into stylized 3D hair cards with colors and spring-bone chains for animation, exported as GLB meshes or a VRM avatar. For avatar make…

VoxCPM2

openbmb · Kit

Speaks text in 30 languages and nine Chinese dialects at 48 kHz, designs a new voice from a written description and clones a voice from a short clip. An open 2…

Diffusion-Qwen3.8-27B

summerAIO · Kit

Generates text by parallel unmasking from Qwen3.8-27B converted into a masked diffusion language model. An experimental research checkpoint for researchers and…

pjit-21k

Linum-AI · Kit

Linum's experimental 2B text to image model, trained for about $21K of compute, that works best with dense 50 to 80 word prompts and cannot write text in image…

MuseTalkDemo

TMElyralab · Kit

A Docker-packaged demo app from TMElyralab for its MuseTalk model, bundled with DWPose, a Stable Diffusion autoencoder and ffmpeg tooling. It suits developers…

ModelScope Text To Video

shamedrank10 · Kit

Turns an English text prompt into a short 256x256 video clip with the ModelScope text-to-video diffusion model, in a Gradio demo. For researchers and curious m…

JoyAI Video Edit

hugging-apps · Kit

Edits a video from a plain-language instruction and streams the result back chunk by chunk while it is still being made. A demo of the JoyAI-Video-Edit model f…

Wet-Ink

Freedomform · Kit

Animates a decoding method for diffusion language models in your browser, where you change a window size and watch tokens stay editable or freeze. For research…

Diffusers Image Outpaint

fffiloni · Kit

Extends a picture past its edges by filling in new content around it, in a Gradio app built on Stable Diffusion XL models. It suits anyone who needs more canva…

LGM-Mini

dylanebert · Kit

Runs a small Gradio demo built on the LGM model and several multi-view diffusion models, on a Hugging Face GPU. It suits people interested in 3D generation res…

Di♪♪Rhythm

ASLP-lab · Kit

Generates songs with the DiffRhythm model from the ASLP lab, in a Gradio Space on GPU hardware. For musicians, researchers and curious users who want to hear A…

Texture-Craft

ArshPathan · Kit

Makes transparent 16x16 item textures from text prompts with a diffusion model, through a Gradio Space, a local REST server or an MCP tool. For game makers and…

SeedVR2 3B Image Upscaler

Upsampler · Kit

Restores and upscales a single image in one diffusion step with the SeedVR2-3B model, in the browser on Hugging Face or from a local copy. Suits anyone with a…

Votepurchase Multiple Model (SD1.5/SDXL Text-to-Image)

John6666 · Kit

Generates images from text prompts with a choice of Stable Diffusion 1.5 and SDXL models on one browser page. It suits people who want to try several image mod…

InstantID

InstantX · Kit

A Gradio demo built around the InstantX/InstantID model, with a Stable Diffusion XL pipeline, ControlNet depth and canny models and style templates, online or…

P-JiT 21K

hugging-apps · Kit

Generates images from text prompts with P-JiT 21K, a diffusion transformer that works directly on pixels. Includes a prompt expander for short prompts. For ima…

Tile Upscaler

gokaygokay · Kit

Enlarges an image with Real-ESRGAN and a diffusion model guided by ControlNet, in a Gradio app. It suits anyone who needs a bigger version of a picture.

RelayVSR Video Upscaler

Upsampler · Kit

Raises the resolution of video with the RelayVSR diffusion model in a Gradio page on Hugging Face. For editors and developers who want sharper footage with tex…

Serverless ImgGen Hub

Nymbo · Kit

Generates images from text prompts with Flux and Stable Diffusion 3.5 models plus style LoRAs, with no GPU needed on your side. Easy to modify. Suits artists,…

Tessera Pixel Studio

thejagstudio · Kit

Paints 16 by 16 pixel art from a text prompt and lets you watch the sprite appear step by step, then edit it. Built around a small diffusion model. Free for pe…